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|---|---|---|---|---|---|---|---|---|---|---|---|---|---|---|---|---|---|---|---|---|---|---|---|---|---|
huggingface/diffusers | 1,272,010,353 | I_kwDOHa8MBc5L0VZx | 10 | https://github.com/huggingface/diffusers/issues/10 | https://api.github.com/repos/huggingface/diffusers/issues/10 | Missing `transformers` requirement | Dependencies don't mention `transformers`
https://github.com/huggingface/diffusers/blob/850d43450f2154d586c1cb226ca17ec53765e2a2/setup.py#L78-L91
However, importing from the quickstart fails without `transformers` installed:
```py
In [1]: import torch
In [2]: from diffusers import UNetModel, DDPMScheduler
... | closed | completed | false | 2 | [] | [
"anton-l"
] | 2022-06-15T10:21:14Z | 2022-06-22T11:48:28Z | 2022-06-22T11:48:28Z | MEMBER | null | 20260407T133413Z | 2026-04-07T13:34:13Z | LysandreJik | 30,755,778 | MDQ6VXNlcjMwNzU1Nzc4 | User | false |
huggingface/diffusers | 1,424,414,610 | I_kwDOHa8MBc5U5teS | 1,000 | https://github.com/huggingface/diffusers/issues/1000 | https://api.github.com/repos/huggingface/diffusers/issues/1000 | `StableDiffusionInpaintPipeline` should handle batches | I'd like to batch inferences with `StableDiffusionInpaintPipeline`. Most of the code in `pipeline_stable_diffusion_inpaint.py` can do it **except** for `prepare_mask_and_masked_image(image, mask)` which **requires** a `PIL.Image`.
**Describe the solution you'd like**
A simple fix would be to just let `tensors` alon... | closed | completed | false | 2 | [
"stale"
] | [] | 2022-10-26T17:17:48Z | 2022-11-30T12:04:46Z | 2022-11-30T12:04:46Z | CONTRIBUTOR | null | 20260407T133413Z | 2026-04-07T13:34:13Z | vict0rsch | 9,283,470 | MDQ6VXNlcjkyODM0NzA= | User | false |
huggingface/diffusers | 2,686,122,280 | I_kwDOHa8MBc6gGvko | 10,003 | https://github.com/huggingface/diffusers/issues/10003 | https://api.github.com/repos/huggingface/diffusers/issues/10003 | DDIMScheduler (or DDPM) of StableDiffusion with 1000 num_inference_steps bugs out | ### Describe the bug
There's something going on with the set_timesteps offset parameters on Stable Diffusion (1v4). The timesteps are set from 1->1000 instead from 0, and so it tries to index out of bounds
### Reproduction
```
model_id = "CompVis/stable-diffusion-v1-4"
scheduler = DDIMScheduler.from_pretrained... | open | null | false | 10 | [
"bug",
"stale"
] | [] | 2024-11-23T15:18:23Z | 2025-02-03T15:04:00Z | null | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | HilaManor | 53,814,901 | MDQ6VXNlcjUzODE0OTAx | User | false |
huggingface/diffusers | 2,686,596,889 | I_kwDOHa8MBc6gIjcZ | 10,004 | https://github.com/huggingface/diffusers/issues/10004 | https://api.github.com/repos/huggingface/diffusers/issues/10004 | how to use kohya sd-scripts flux loras with text encoder keys in diffusers? | resulting lora weights from setting train text encoder to true is incompatible with diffusers load_lora_weights. the script networks/convert_flux_lora.py does not convert the text encoder keys either. | open | null | false | 12 | [
"contributions-welcome"
] | [] | 2024-11-23T20:54:30Z | 2025-03-16T15:39:25Z | null | CONTRIBUTOR | null | 20260407T133413Z | 2026-04-07T13:34:13Z | neuron-party | 96,799,331 | U_kgDOBcUKYw | User | false |
huggingface/diffusers | 2,689,130,458 | I_kwDOHa8MBc6gSN_a | 10,010 | https://github.com/huggingface/diffusers/issues/10010 | https://api.github.com/repos/huggingface/diffusers/issues/10010 | FLUX Fill example not working due to Not Found for url: https://huggingface.co/black-forest-labs/FLUX.1-Fill-dev/resolve/main/model_index.json | ### Describe the bug
FLUX Fill is example not working due to Not Found for url: https://huggingface.co/black-forest-labs/FLUX.1-Fill-dev/resolve/main/model_index.json
### Reproduction
Execute the FLUX Fill example according to doc.
### Logs
_No response_
### System Info
diffusers main
### Who can help?
_No res... | closed | completed | false | 2 | [
"bug"
] | [] | 2024-11-25T04:50:06Z | 2024-11-25T04:53:33Z | 2024-11-25T04:53:33Z | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | petergerten | 17,756,065 | MDQ6VXNlcjE3NzU2MDY1 | User | false |
huggingface/diffusers | 2,690,235,944 | I_kwDOHa8MBc6gWb4o | 10,013 | https://github.com/huggingface/diffusers/issues/10013 | https://api.github.com/repos/huggingface/diffusers/issues/10013 | [Core] refactor model loading | Currently, we have got two codepaths:
1. For non-sharded checkpoints we do: https://github.com/huggingface/diffusers/blob/047bf492914ddc9393070b8f73bba5ad5823eb29/src/diffusers/models/modeling_utils.py#L855
2. For sharded checkpoints we do: https://github.com/huggingface/diffusers/blob/047bf492914ddc9393070b8f73bba... | closed | completed | false | 2 | [
"wip",
"roadmap"
] | [
"sayakpaul",
"SunMarc"
] | 2024-11-25T10:44:58Z | 2025-02-19T12:04:55Z | 2025-02-19T12:04:55Z | MEMBER | null | 20260407T133413Z | 2026-04-07T13:34:13Z | sayakpaul | 22,957,388 | MDQ6VXNlcjIyOTU3Mzg4 | User | false |
huggingface/diffusers | 2,690,920,364 | I_kwDOHa8MBc6gZC-s | 10,016 | https://github.com/huggingface/diffusers/issues/10016 | https://api.github.com/repos/huggingface/diffusers/issues/10016 | Cannot load SD3.5M with from_single_file, mismatched shape for pos_embed.pos_embed in empty state dict | ### Describe the bug
When trying to load the Stable Diffusion 3.5 Medium checkpoint from Stability AI with the following command:
pipe = diffusers.StableDiffusion3Pipeline.from_single_file( '...path/sd3.5_medium.safetensors', text_encoder=None, text_encoder_2=None, text_encoder_3=None )
The following error is ... | closed | completed | false | 3 | [
"bug"
] | [
"DN6"
] | 2024-11-25T14:25:18Z | 2024-12-11T07:41:41Z | 2024-12-03T01:21:42Z | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | CodeExplode | 79,982,924 | MDQ6VXNlcjc5OTgyOTI0 | User | false |
huggingface/diffusers | 2,691,321,209 | I_kwDOHa8MBc6gak15 | 10,017 | https://github.com/huggingface/diffusers/issues/10017 | https://api.github.com/repos/huggingface/diffusers/issues/10017 | loading dataset issue with load_dataset() | ### Describe the bug
i'm unable to load my dataset for [controlnet training](https://github.com/huggingface/diffusers/blob/074e12358bc17e7dbe111ea4f62f05dbae8a49d5/examples/controlnet/train_controlnet.py#L606) using load_dataset(). however, load_from_disk() seems to work?
would appreciate if someone can explain wh... | closed | completed | false | 1 | [
"bug"
] | [] | 2024-11-25T16:15:45Z | 2024-11-26T03:32:51Z | 2024-11-26T03:32:51Z | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | sarahahtee | 81,594,044 | MDQ6VXNlcjgxNTk0MDQ0 | User | false |
huggingface/diffusers | 2,691,771,169 | I_kwDOHa8MBc6gcSsh | 10,018 | https://github.com/huggingface/diffusers/issues/10018 | https://api.github.com/repos/huggingface/diffusers/issues/10018 | no model_index.json for flux tools | cannot load any flux tools pipelines (fluxfill, fluxcontrol, etc) all throw a model_index.json missing error. | closed | completed | false | 2 | [] | [] | 2024-11-25T18:22:53Z | 2024-11-25T23:55:23Z | 2024-11-25T23:55:22Z | CONTRIBUTOR | null | 20260407T133413Z | 2026-04-07T13:34:13Z | neuron-party | 96,799,331 | U_kgDOBcUKYw | User | false |
huggingface/diffusers | 2,692,979,032 | I_kwDOHa8MBc6gg5lY | 10,022 | https://github.com/huggingface/diffusers/issues/10022 | https://api.github.com/repos/huggingface/diffusers/issues/10022 | [core] refactor `attention_processor.py` the easy way | With @DN6 we have been discussing an idea about breaking up `src/diffusers/models/attention_processor.py` as it's getting excruciatingly longer and longer. The idea is simple and won't very likely require multiple rounds of PRs.
* Create a module named `attention_processor`.
* Split the attention processor classes... | open | null | false | 6 | [
"wip"
] | [
"DN6"
] | 2024-11-26T03:37:33Z | 2026-04-01T07:55:23Z | null | MEMBER | null | 20260407T133413Z | 2026-04-07T13:34:13Z | sayakpaul | 22,957,388 | MDQ6VXNlcjIyOTU3Mzg4 | User | false |
huggingface/diffusers | 2,692,982,333 | I_kwDOHa8MBc6gg6Y9 | 10,023 | https://github.com/huggingface/diffusers/issues/10023 | https://api.github.com/repos/huggingface/diffusers/issues/10023 | Parallel Inference support ffor Dit model | does diffusers have any plan to support Parallel Inference like TP/PP/SP? | open | null | false | 3 | [
"stale"
] | [] | 2024-11-26T03:40:48Z | 2024-12-26T15:02:57Z | null | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | enozhu | 37,290,567 | MDQ6VXNlcjM3MjkwNTY3 | User | false |
huggingface/diffusers | 2,693,777,714 | I_kwDOHa8MBc6gj8ky | 10,025 | https://github.com/huggingface/diffusers/issues/10025 | https://api.github.com/repos/huggingface/diffusers/issues/10025 | attention mask for transformer Flux | ### Describe the bug
Is it possible to get back the `attention_mask` argument in the flux attention processor
```
hidden_states = F.scaled_dot_product_attention(query, key, value, dropout_p=0.0, is_causal=False,attn_mask=attention_mask)
```
https://github.com/huggingface/diffusers/blob/main/src/diffusers/mo... | closed | completed | false | 19 | [
"bug"
] | [] | 2024-11-26T08:51:20Z | 2024-12-05T00:22:37Z | 2024-12-05T00:22:37Z | CONTRIBUTOR | null | 20260407T133413Z | 2026-04-07T13:34:13Z | christopher5106 | 6,875,375 | MDQ6VXNlcjY4NzUzNzU= | User | false |
huggingface/diffusers | 2,693,996,392 | I_kwDOHa8MBc6gkx9o | 10,026 | https://github.com/huggingface/diffusers/issues/10026 | https://api.github.com/repos/huggingface/diffusers/issues/10026 | Flux ControlNet Training Multi-GPU DeepSpeed Stage-3 doesn't reduce memory compared to Single GPU | ### Describe the bug
I am running a slightly modified version of [Flux ControlNet training script in diffusers](https://github.com/huggingface/diffusers/blob/main/examples/controlnet/train_controlnet_flux.py). The script is attached below. I am using DeepSpeed Stage-3 with the accelerate config below.
When I use on... | closed | completed | true | 1 | [
"bug"
] | [] | 2024-11-26T09:59:56Z | 2024-11-26T11:13:26Z | 2024-11-26T11:13:26Z | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | enesmsahin | 36,894,859 | MDQ6VXNlcjM2ODk0ODU5 | User | false |
huggingface/diffusers | 2,695,890,786 | I_kwDOHa8MBc6gsAdi | 10,029 | https://github.com/huggingface/diffusers/issues/10029 | https://api.github.com/repos/huggingface/diffusers/issues/10029 | Linear' object has no attribute 'base_layer | ### Describe the bug
I was trying to inference flux lora , lora was finetuned on ai toolkit
I have installed diffusers from source
while inferencing with lora , I am getting this error ===> Linear' object has no attribute 'base_layer
Please let me know if you need any info
Thanks
... | open | null | false | 2 | [
"bug",
"stale"
] | [] | 2024-11-26T19:37:55Z | 2024-12-27T15:02:59Z | null | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | Sandeep-Narahari | 46,187,761 | MDQ6VXNlcjQ2MTg3NzYx | User | false |
huggingface/diffusers | 2,696,254,448 | I_kwDOHa8MBc6gtZPw | 10,030 | https://github.com/huggingface/diffusers/issues/10030 | https://api.github.com/repos/huggingface/diffusers/issues/10030 | ControlNet expects added_cond_kwargs to not be null, which pipeline_controlnet_img2img does not respect | ### Describe the bug
Please see this code in pipeline_control_img2img.py
https://github.com/huggingface/diffusers/blob/ad5ecd1251472dbc69da1268671d41bc2d8c1caa/src/diffusers/pipelines/controlnet/pipeline_controlnet_img2img.py#L1199-L1203
### Reproduction
I think the pipeline is just straight up bugged sin... | closed | completed | false | 3 | [
"bug"
] | [] | 2024-11-26T22:19:13Z | 2024-11-26T22:55:35Z | 2024-11-26T22:55:35Z | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | oxysoft | 3,945,277 | MDQ6VXNlcjM5NDUyNzc= | User | false |
huggingface/diffusers | 2,697,762,971 | I_kwDOHa8MBc6gzJib | 10,034 | https://github.com/huggingface/diffusers/issues/10034 | https://api.github.com/repos/huggingface/diffusers/issues/10034 | Add Support for Lightricks/LTX-Video Model and Video Generation Features (text2video, image2video, video2video) in Diffusers with Sequential CPU Offload Compatibility | **Is your feature request related to a problem? Please describe.**
Currently, the Diffusers library does not provide native support for advanced video generation tasks such as text-to-video (text2video), image-to-video (image2video), or video-to-video (video2video) generation for lighttricks/LTX-Video model. With the ... | closed | completed | false | 2 | [] | [] | 2024-11-27T09:10:41Z | 2024-11-28T05:33:37Z | 2024-11-28T05:33:37Z | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | BasimBashir | 45,428,008 | MDQ6VXNlcjQ1NDI4MDA4 | User | false |
huggingface/diffusers | 2,698,659,585 | I_kwDOHa8MBc6g2kcB | 10,035 | https://github.com/huggingface/diffusers/issues/10035 | https://api.github.com/repos/huggingface/diffusers/issues/10035 | CogvideoX VAE decoder consumes significantly more memory in the latest version | ### Describe the bug
The memory consumption for CogVideoX decoder in diffusers 0.31.0 version consumes significantly more memory. To the point where the model goes OOM even on 80G H100 GPUs with a relatively modest frame count.
I include two profiles for very small input tensors of only 5 frames where its visible ... | closed | completed | true | 1 | [
"bug"
] | [] | 2024-11-27T14:02:53Z | 2024-11-27T15:49:18Z | 2024-11-27T15:49:18Z | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | ic-synth | 170,024,644 | U_kgDOCiJexA | User | false |
huggingface/diffusers | 2,701,086,646 | I_kwDOHa8MBc6g_0-2 | 10,039 | https://github.com/huggingface/diffusers/issues/10039 | https://api.github.com/repos/huggingface/diffusers/issues/10039 | Forgot to pass the output sample size in AttnUpBlock2D() in unet_2d_blocks.py | ### Describe the bug
Hi, I found that the parameter upsample_size was not passed into the upsampler in AttnUpBlock2D() in unet_2d_blocks.py [Line: 2386]. This can result in mismatched dimensions when processing images using AttnUpBlock2D(). This does not seem to have been corrected in the latest GitHub version [v0.31.... | closed | completed | false | 3 | [
"bug",
"stale"
] | [] | 2024-11-28T07:14:19Z | 2024-12-28T19:58:17Z | 2024-12-28T19:58:16Z | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | xyliu6666 | 92,036,814 | U_kgDOBXxezg | User | false |
huggingface/diffusers | 2,701,283,255 | I_kwDOHa8MBc6hAk-3 | 10,040 | https://github.com/huggingface/diffusers/issues/10040 | https://api.github.com/repos/huggingface/diffusers/issues/10040 | [BUG - STABLE DIFFUSION 3] Grey images generated | ### Describe the bug
I'm running the SD3 model [stabilityai/stable-diffusion-3-medium](https://huggingface.co/stabilityai/stable-diffusion-3-medium) with the following settings:
- Height: 1024
- Width: 1024
- Inference steps: 50
- Guidance scale: 7
- Prompts length: 32 (using the same input prompt: "A men jumps... | closed | completed | false | 13 | [
"bug"
] | [] | 2024-11-28T08:31:27Z | 2024-12-02T23:33:58Z | 2024-12-02T23:33:58Z | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | hoangvictor | 41,613,226 | MDQ6VXNlcjQxNjEzMjI2 | User | false |
huggingface/diffusers | 2,701,318,031 | I_kwDOHa8MBc6hAteP | 10,041 | https://github.com/huggingface/diffusers/issues/10041 | https://api.github.com/repos/huggingface/diffusers/issues/10041 | 🙋Flexible Directory Structure Support for Custom Pipeline Scenarios | ## 🚨 Flexible Directory Structure for Custom Pipeline Scenarios Raised Error!
When I have defined some custom_pipelines in local and the directory structure looks like:
```bash
custom_pipelines
├── custom_modules
│ ├── autoencoders
│ │ └── vae_a.py
│ └── transformers
│ ├── transformer_a.py
│ ... | open | null | false | 4 | [
"stale"
] | [
"a-r-r-o-w"
] | 2024-11-28T08:45:20Z | 2025-01-22T15:03:43Z | null | CONTRIBUTOR | null | 20260407T133413Z | 2026-04-07T13:34:13Z | townwish4git | 143,256,262 | U_kgDOCInqxg | User | false |
huggingface/diffusers | 2,701,802,452 | I_kwDOHa8MBc6hCjvU | 10,043 | https://github.com/huggingface/diffusers/issues/10043 | https://api.github.com/repos/huggingface/diffusers/issues/10043 | F5-TTS Integration | ### Model/Pipeline/Scheduler description
F5-TTS is a fully non-autoregressive text-to-speech system based on flow matching with Diffusion Transformer (DiT).
It has excellent voice cloning capabilities, and audio generation is of quite high quality.
### Open source status
- [X] The model implementation is available.... | open | reopened | false | 11 | [
"help wanted",
"contributions-welcome"
] | [] | 2024-11-28T11:14:18Z | 2025-11-02T18:46:02Z | null | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | nityanandmathur | 77,379,835 | MDQ6VXNlcjc3Mzc5ODM1 | User | false |
huggingface/diffusers | 2,702,055,025 | I_kwDOHa8MBc6hDhZx | 10,045 | https://github.com/huggingface/diffusers/issues/10045 | https://api.github.com/repos/huggingface/diffusers/issues/10045 | Problem with importing pipelines | ### Describe the bug
When trying to import pipelines in colab i have the following error:
RuntimeError Traceback (most recent call last)
[/usr/local/lib/python3.10/dist-packages/transformers/utils/import_utils.py](https://localhost:8080/#) in _get_module(self, module_name)
1777 ... | closed | completed | false | 0 | [
"bug"
] | [] | 2024-11-28T12:50:50Z | 2024-11-28T18:03:57Z | 2024-11-28T18:03:57Z | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | stickman712 | 159,720,673 | U_kgDOCYUk4Q | User | false |
huggingface/diffusers | 2,702,916,663 | I_kwDOHa8MBc6hGzw3 | 10,049 | https://github.com/huggingface/diffusers/issues/10049 | https://api.github.com/repos/huggingface/diffusers/issues/10049 | Cannot import MochiPipeline | ### Describe the bug
Can't import mochi in colab. When trying have a following error:
ImportError Traceback (most recent call last)
[<ipython-input-5-6aeec58dc926>](https://localhost:8080/#) in <cell line: 2>()
1 import torch
----> 2 from diffusers import MochiPipeline
... | closed | completed | false | 1 | [
"bug"
] | [] | 2024-11-28T18:11:49Z | 2024-11-28T18:21:12Z | 2024-11-28T18:21:12Z | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | stickman712 | 159,720,673 | U_kgDOCYUk4Q | User | false |
huggingface/diffusers | 1,424,800,551 | I_kwDOHa8MBc5U7Lsn | 1,005 | https://github.com/huggingface/diffusers/issues/1005 | https://api.github.com/repos/huggingface/diffusers/issues/1005 | num_inference_timesteps can't exceed num_train_timesteps? | ### Describe the bug
Hey all,
I recently trained an unconditional image generation DDPM pipeline on a custom dataset using the `train_unconditional.py` script. The scheduler specified during training had `num_train_timesteps=1000` as follows:
```
noise_scheduler = DDPMScheduler(num_train_timesteps=1000)
```
... | closed | completed | false | 2 | [
"bug"
] | [] | 2022-10-26T23:08:36Z | 2022-10-30T02:43:50Z | 2022-10-30T02:43:50Z | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | mihail911 | 2,789,441 | MDQ6VXNlcjI3ODk0NDE= | User | false |
huggingface/diffusers | 2,704,257,129 | I_kwDOHa8MBc6hL7Bp | 10,050 | https://github.com/huggingface/diffusers/issues/10050 | https://api.github.com/repos/huggingface/diffusers/issues/10050 | Is there any img2img KDiffusion equivalent of StableDiffusionKDiffusionPipeline? | ### Model/Pipeline/Scheduler description
I'm working on result alignment between diffusers and A1111 webui.
In txt2img scene, I can achieve via `StableDiffusionKDiffusionPipeline`, refer to https://github.com/huggingface/diffusers/issues/3253.
But in img2img scene, is there any KDiffusion pipeline equivalent?
I... | open | null | false | 2 | [
"stale"
] | [] | 2024-11-29T07:47:11Z | 2024-12-29T15:03:05Z | null | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | juju812 | 19,584,910 | MDQ6VXNlcjE5NTg0OTEw | User | false |
huggingface/diffusers | 2,704,631,827 | I_kwDOHa8MBc6hNWgT | 10,052 | https://github.com/huggingface/diffusers/issues/10052 | https://api.github.com/repos/huggingface/diffusers/issues/10052 | syntax error on ctx | ### Describe the bug
Across [several files](https://github.com/search?q=repo%3Ahuggingface%2Fdiffusers%20CTX&type=code), a bool statement is wrongly constructed IMO
One example: https://github.com/huggingface/diffusers/blob/c96bfa5c80eca798d555a79a491043c311d0f608/scripts/convert_sd3_to_diffusers.py#L14 => instea... | closed | completed | false | 2 | [
"bug"
] | [] | 2024-11-29T10:13:59Z | 2024-12-02T05:58:02Z | 2024-12-02T05:58:01Z | CONTRIBUTOR | null | 20260407T133413Z | 2026-04-07T13:34:13Z | Clement-Lelievre | 70,368,164 | MDQ6VXNlcjcwMzY4MTY0 | User | false |
huggingface/diffusers | 2,704,908,993 | I_kwDOHa8MBc6hOaLB | 10,054 | https://github.com/huggingface/diffusers/issues/10054 | https://api.github.com/repos/huggingface/diffusers/issues/10054 | Testing beta sigmas in EulerDiscreteScheduler passes wrong parameters to torch.linspace | ### Describe the bug
Creating a EulerDiscreteScheduler using an existing schedulers config seems to pass a `FrozenDict` for `num_train_timesteps` instead of the expected `int`
e.g.
```py
pipe.scheduler = EulerDiscreteScheduler(pipe.scheduler.config, timestep_spacing="linspace", use_beta_sigmas=True)
```
fails w... | closed | completed | false | 1 | [
"bug"
] | [] | 2024-11-29T11:47:28Z | 2024-11-29T11:55:34Z | 2024-11-29T11:55:33Z | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | Vargol | 62,868 | MDQ6VXNlcjYyODY4 | User | false |
huggingface/diffusers | 2,705,183,952 | I_kwDOHa8MBc6hPdTQ | 10,055 | https://github.com/huggingface/diffusers/issues/10055 | https://api.github.com/repos/huggingface/diffusers/issues/10055 | Training script for a Controlnet based on SD3 does not work | ### Describe the bug
Hi @sayakpaul and all others :)
The training script for a Control-net based on Stable Diffusion 3 seems to not work.
**RuntimeError: Given groups=1, weight of size [1536, 17, 2, 2], expected input[4, 16, 64, 64] to have 17 channels, but got 16 channels instead**
I tried to follow th... | open | null | false | 17 | [
"bug",
"stale"
] | [] | 2024-11-29T13:46:29Z | 2025-02-03T15:03:46Z | null | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | Putzzmunta | 29,607,612 | MDQ6VXNlcjI5NjA3NjEy | User | false |
huggingface/diffusers | 2,705,420,279 | I_kwDOHa8MBc6hQW_3 | 10,057 | https://github.com/huggingface/diffusers/issues/10057 | https://api.github.com/repos/huggingface/diffusers/issues/10057 | Token Issue with FluxTransformer2DModel.from_single_file Method | ### Describe the bug
When I use the FluxTransformer2DModel.from_single_file method, the provided token seems to be invalid, and I am unable to download files that require token verification, such as the official Flux model files.
However, I can confirm that the token is valid because when using FluxPipeline.from_... | closed | completed | false | 8 | [
"bug"
] | [] | 2024-11-29T15:18:06Z | 2024-12-28T17:06:37Z | 2024-12-03T01:15:34Z | CONTRIBUTOR | null | 20260407T133413Z | 2026-04-07T13:34:13Z | CyberVy | 72,680,847 | MDQ6VXNlcjcyNjgwODQ3 | User | false |
huggingface/diffusers | 2,706,472,451 | I_kwDOHa8MBc6hUX4D | 10,059 | https://github.com/huggingface/diffusers/issues/10059 | https://api.github.com/repos/huggingface/diffusers/issues/10059 | [core] introduce `AutoModel` | Now that we have quite a few model classes (and transformer model classes), I wonder if it makes sense to have an AutoModel class. Oftentimes, users initialize the model separately and pass it to a `DiffusionPipeline`:
```py
transformer = "..."
pipeline = DiffusionPipeline.from_pretrained("black-forest-labs/FLUX.1... | closed | completed | false | 0 | [
"enhancement"
] | [] | 2024-11-30T01:24:39Z | 2025-04-09T09:50:10Z | 2025-04-09T09:50:10Z | MEMBER | null | 20260407T133413Z | 2026-04-07T13:34:13Z | sayakpaul | 22,957,388 | MDQ6VXNlcjIyOTU3Mzg4 | User | false |
huggingface/diffusers | 2,707,499,364 | I_kwDOHa8MBc6hYSlk | 10,062 | https://github.com/huggingface/diffusers/issues/10062 | https://api.github.com/repos/huggingface/diffusers/issues/10062 | 'UNet2DConditionModel' object has no attribute 'peft_config' | ### Describe the bug
After pulling the latest code from the _examples/research_projects/instructpix2pix_lora_ repository, I encountered the following error when running the default code:
`AttributeError: 'Linear' object has no attribute 'set_lora_layer'`
I resolved this issue by modifying the code.
The training now... | open | null | false | 6 | [
"bug",
"stale"
] | [
"Aiden-Frost"
] | 2024-11-30T15:25:58Z | 2025-01-06T15:03:26Z | null | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | duduke121 | 144,907,684 | U_kgDOCKMdpA | User | false |
huggingface/diffusers | 2,708,432,047 | I_kwDOHa8MBc6hb2Sv | 10,066 | https://github.com/huggingface/diffusers/issues/10066 | https://api.github.com/repos/huggingface/diffusers/issues/10066 | self.attn ValueError: not enough values to unpack (expected 2, got 1) | ### Describe the bug
I have encountered an error when trying to sample from SD3 transformer component, where line 208 in `attention.py` is asking for two outputs but only one is given:
```
File "/home/p2p/pytorch/mc3/envs/marigold/lib/python3.10/site-packages/diffusers/src/diffusers/models/attention.py", line ... | open | null | false | 15 | [
"bug",
"stale",
"needs-code-example"
] | [] | 2024-12-01T03:59:35Z | 2025-06-19T07:23:49Z | null | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | tracyyy-zhu | 61,514,951 | MDQ6VXNlcjYxNTE0OTUx | User | false |
huggingface/diffusers | 2,709,432,084 | I_kwDOHa8MBc6hfqcU | 10,070 | https://github.com/huggingface/diffusers/issues/10070 | https://api.github.com/repos/huggingface/diffusers/issues/10070 | Xformers info , memory efficient atttention unavailable | ### Describe the bug
I just started learning Stable Diffuision on Win11. After I installed xformers, I found several memory_efficient_attention string is unavailable. Is it possible to make them available? Thanks for any help.
### Reproduction
xFormers 0.0.28.post3
memory_efficient_attention.ckF: ... | open | null | false | 1 | [
"bug",
"stale"
] | [] | 2024-12-01T16:14:21Z | 2025-01-01T15:03:09Z | null | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | Stareshine | 160,834,677 | U_kgDOCZYkdQ | User | false |
huggingface/diffusers | 2,709,616,761 | I_kwDOHa8MBc6hgXh5 | 10,071 | https://github.com/huggingface/diffusers/issues/10071 | https://api.github.com/repos/huggingface/diffusers/issues/10071 | [New feature] A Noise Injection Method for Flux | ### New Feature for FluxPipeline
Paper: [Link](https://openreview.net/pdf?id=KzokzKV4JK)
Code : [Link](https://anonymous.4open.science/r/SSNI-F746/README.md)
workflow : [Link](https://www.dropbox.com/scl/fi/hhitjx6lqpqpv8xjx9ikq/Flux-Noise-Injection.json?rlkey=45xnu45j1i5owiwhc7z1hppn1&e=1&dl=0)
This paper introduc... | open | null | false | 10 | [
"wip",
"consider-for-modular-diffusers"
] | [] | 2024-12-01T18:35:39Z | 2025-02-26T16:36:07Z | null | CONTRIBUTOR | null | 20260407T133413Z | 2026-04-07T13:34:13Z | SahilCarterr | 110,806,554 | U_kgDOBprGGg | User | false |
huggingface/diffusers | 2,710,794,006 | I_kwDOHa8MBc6hk28W | 10,074 | https://github.com/huggingface/diffusers/issues/10074 | https://api.github.com/repos/huggingface/diffusers/issues/10074 | how to install diffusers 0.32.0 | FluxFillPipeline Function need =0.32.0 But I don't know how to install it, can anyone help me? Thanks in advance | closed | completed | false | 2 | [] | [] | 2024-12-02T07:05:24Z | 2024-12-02T19:11:34Z | 2024-12-02T09:35:05Z | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | babyta | 37,787,795 | MDQ6VXNlcjM3Nzg3Nzk1 | User | false |
huggingface/diffusers | 2,710,824,259 | I_kwDOHa8MBc6hk-VD | 10,075 | https://github.com/huggingface/diffusers/issues/10075 | https://api.github.com/repos/huggingface/diffusers/issues/10075 | flux:different results on different machines. | Hi!
During the debugging phase, I've noticed that training the [FLUX-Controlnet-Inpainting](https://github.com/alimama-creative/FLUX-Controlnet-Inpainting) model yields different results on different machines.
Specifically, on one machine, everything trains fine,
but on another, I'm getting an error that says "... | open | null | false | 3 | [
"stale"
] | [] | 2024-12-02T07:22:00Z | 2025-01-01T15:03:03Z | null | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | D222097 | 50,061,868 | MDQ6VXNlcjUwMDYxODY4 | User | false |
huggingface/diffusers | 2,710,923,544 | I_kwDOHa8MBc6hlWkY | 10,076 | https://github.com/huggingface/diffusers/issues/10076 | https://api.github.com/repos/huggingface/diffusers/issues/10076 | Do we have any script covert from hf format to orginal format? | **Is your feature request related to a problem? Please describe.**
scripts/convert_cogvideox_to_diffusers.py
in this script, we can convert cogvideox -> diffusers. Do we have the opposite script?
cc @yiyixuxu
| open | null | false | 1 | [
"good first issue",
"contributions-welcome",
"conversion script"
] | [] | 2024-12-02T07:49:34Z | 2024-12-02T18:22:50Z | null | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | foreverpiano | 29,483,524 | MDQ6VXNlcjI5NDgzNTI0 | User | false |
huggingface/diffusers | 2,711,410,316 | I_kwDOHa8MBc6hnNaM | 10,078 | https://github.com/huggingface/diffusers/issues/10078 | https://api.github.com/repos/huggingface/diffusers/issues/10078 | Cache text encoder embeds in pipelines | **Is your feature request related to a problem? Please describe.**
When reusing a prompt text encoder embeds are recomputed, this can be time consuming for something like T5-XXL with offloading or on CPU.
Text encoder embeds are relatively small, so keeping them in memory is feasible.
```python
import torch
... | open | null | false | 3 | [
"community-examples",
"stale"
] | [] | 2024-12-02T10:40:01Z | 2025-01-01T15:02:59Z | null | CONTRIBUTOR | null | 20260407T133413Z | 2026-04-07T13:34:13Z | hlky | 106,811,348 | U_kgDOBl3P1A | User | false |
huggingface/diffusers | 2,711,545,173 | I_kwDOHa8MBc6hnuVV | 10,080 | https://github.com/huggingface/diffusers/issues/10080 | https://api.github.com/repos/huggingface/diffusers/issues/10080 | Proposal to add sigmas option to FluxPipeline | **Is your feature request related to a problem? Please describe.**
Flux's image generation is great, but it seems to have a tendency to remove too much detail.
**Describe the solution you'd like.**
Add the `sigmas` option to `FluxPipeline` to enable adjustment of the degree of noise removal.
**Additional contex... | closed | completed | false | 3 | [] | [] | 2024-12-02T11:14:40Z | 2024-12-02T18:16:49Z | 2024-12-02T18:16:49Z | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | John6666cat | 186,692,226 | U_kgDOCyCygg | User | false |
huggingface/diffusers | 2,711,831,082 | I_kwDOHa8MBc6ho0Iq | 10,085 | https://github.com/huggingface/diffusers/issues/10085 | https://api.github.com/repos/huggingface/diffusers/issues/10085 | RuntimeError with LyCORIS, Batch Inference and skip_guidance_layers | ### Describe the bug
A `RuntimeError` occurs when using the following combination:
* SD3
* Batch inference (`num_images_per_prompt > 1`)
* LyCORIS
* `skip_guidance_layers` is set
The error message is: `"RuntimeError: The size of tensor a (2) must match the size of tensor b (4) at non-singleton dimension 0... | closed | completed | false | 5 | [
"bug"
] | [] | 2024-12-02T12:56:30Z | 2024-12-02T21:43:05Z | 2024-12-02T21:43:04Z | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | japppie | 131,072,659 | U_kgDOB9ACkw | User | false |
huggingface/diffusers | 1,424,872,662 | I_kwDOHa8MBc5U7dTW | 1,009 | https://github.com/huggingface/diffusers/issues/1009 | https://api.github.com/repos/huggingface/diffusers/issues/1009 | Weight decay bug in textual_inversion.py | ### Describe the bug
I think there's a bug in the [textual_inversion.py ](https://github.com/huggingface/diffusers/blob/main/examples/textual_inversion/textual_inversion.py)script. Specifically, the word embeddings of all the tokens except the newly added token are supposed to be frozen. The way the script does this i... | closed | completed | false | 5 | [
"bug",
"stale"
] | [
"patil-suraj"
] | 2022-10-27T00:57:36Z | 2023-01-13T15:05:01Z | 2023-01-13T15:05:01Z | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | scoutsaachi | 9,390,538 | MDQ6VXNlcjkzOTA1Mzg= | User | false |
huggingface/diffusers | 2,713,832,446 | I_kwDOHa8MBc6hwcv- | 10,090 | https://github.com/huggingface/diffusers/issues/10090 | https://api.github.com/repos/huggingface/diffusers/issues/10090 | Face editor onnxruntime error Fail:Loadlibrary faild with onnxruntime_providers_cuda.dll | ### Describe the bug
After I installed the Face Editor in Webui/Extension. the starter show me such fails. Even though, Face editor will still work fine. That confuses me. These will now show when I disable this extensions.
I am a beginner. I have checked everything I could find on the internet. Onnxruntime-GPU 1.20.... | closed | completed | false | 3 | [
"bug",
"stale"
] | [] | 2024-12-03T02:32:31Z | 2025-01-02T18:50:22Z | 2025-01-02T18:50:22Z | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | Stareshine | 160,834,677 | U_kgDOCZYkdQ | User | false |
huggingface/diffusers | 2,713,868,101 | I_kwDOHa8MBc6hwldF | 10,091 | https://github.com/huggingface/diffusers/issues/10091 | https://api.github.com/repos/huggingface/diffusers/issues/10091 | CogvideoX VAE decoder consumes significantly more memory in the latest version | ### Discussed in https://github.com/huggingface/diffusers/discussions/10036
<div type='discussions-op-text'>
<sup>Originally posted by **ic-synth** November 27, 2024</sup>
### Describe the bug
The memory consumption for CogVideoX decoder in diffusers 0.31.0 version consumes significantly more memory. To the p... | open | null | false | 6 | [
"bug",
"wip",
"roadmap"
] | [
"yiyixuxu",
"a-r-r-o-w"
] | 2024-12-03T03:02:42Z | 2025-08-10T18:29:54Z | null | MEMBER | null | 20260407T133413Z | 2026-04-07T13:34:13Z | sayakpaul | 22,957,388 | MDQ6VXNlcjIyOTU3Mzg4 | User | false |
huggingface/diffusers | 2,713,961,860 | I_kwDOHa8MBc6hw8WE | 10,093 | https://github.com/huggingface/diffusers/issues/10093 | https://api.github.com/repos/huggingface/diffusers/issues/10093 | [pipelines] add better checking when a wrong model is passed when initializing a pipeline | With this code:
```py
from diffusers import DiffusionPipeline
from transformers import T5EncoderModel
import torch
repo_id = "black-forest-labs/FLUX.1-dev"
text_encoder_2 = T5EncoderModel.from_pretrained(repo_id, subfolder="text_encoder_2")
pipe = DiffusionPipeline.from_pretrained(
repo_id, text_enco... | closed | completed | false | 1 | [
"roadmap"
] | [] | 2024-12-03T04:33:35Z | 2024-12-19T09:24:53Z | 2024-12-19T09:24:53Z | MEMBER | null | 20260407T133413Z | 2026-04-07T13:34:13Z | sayakpaul | 22,957,388 | MDQ6VXNlcjIyOTU3Mzg4 | User | false |
huggingface/diffusers | 2,714,263,146 | I_kwDOHa8MBc6hyF5q | 10,097 | https://github.com/huggingface/diffusers/issues/10097 | https://api.github.com/repos/huggingface/diffusers/issues/10097 | Missing Section in Optimization Docs | ### Describe the bug
There is a missing section in the Optimization section of the diffuser docs. Its location being here in the ReadMe:

Web link for the main version: https://huggingface.co/docs/diffusers/main/en/optimi... | closed | completed | false | 2 | [
"bug"
] | [] | 2024-12-03T07:56:20Z | 2024-12-03T18:23:36Z | 2024-12-03T18:23:36Z | CONTRIBUTOR | null | 20260407T133413Z | 2026-04-07T13:34:13Z | ParagEkbote | 69,567,729 | MDQ6VXNlcjY5NTY3NzI5 | User | false |
huggingface/diffusers | 2,714,883,731 | I_kwDOHa8MBc6h0daT | 10,100 | https://github.com/huggingface/diffusers/issues/10100 | https://api.github.com/repos/huggingface/diffusers/issues/10100 | Add support for ConsisID | ### Model/Pipeline/Scheduler description
I have found a awesome work ConsisID, a tuning-free DiT-based controllable IPT2V model to keep human-identity consistent in the generated video. Can diffuser support it, so that I can simply use this model, thanks
### Open source status
- [X] The model implementation is... | closed | completed | false | 1 | [
"New pipeline/model"
] | [] | 2024-12-03T12:21:01Z | 2025-06-10T20:55:00Z | 2025-06-10T20:54:59Z | CONTRIBUTOR | null | 20260407T133413Z | 2026-04-07T13:34:13Z | 963658029 | 25,529,908 | MDQ6VXNlcjI1NTI5OTA4 | User | false |
huggingface/diffusers | 2,715,881,329 | I_kwDOHa8MBc6h4Q9x | 10,106 | https://github.com/huggingface/diffusers/issues/10106 | https://api.github.com/repos/huggingface/diffusers/issues/10106 | Add HunyuanVideo | ## HunyuanVideo
_We present HunyuanVideo, a novel open-source video foundation model that exhibits performance in video generation that is comparable to, if not superior to, leading closed-source models. HunyuanVideo features a comprehensive framework that integrates several key contributions, including data curatio... | closed | completed | false | 12 | [
"New pipeline/model"
] | [] | 2024-12-03T20:09:15Z | 2024-12-31T06:30:17Z | 2024-12-17T21:39:19Z | CONTRIBUTOR | null | 20260407T133413Z | 2026-04-07T13:34:13Z | hlky | 106,811,348 | U_kgDOBl3P1A | User | false |
huggingface/diffusers | 2,716,073,827 | I_kwDOHa8MBc6h4_9j | 10,107 | https://github.com/huggingface/diffusers/issues/10107 | https://api.github.com/repos/huggingface/diffusers/issues/10107 | gradient checkpointing runs during validations in the training examples | ### Describe the bug
The gradient checkpointing is enabled via self.training, but then the log_validations also unnecessarily encounter this codepath.
I found that I have to disable this when running even under `torch.no_grad()`, looked and saw that the official examples do not do this either.
This gives a sub... | closed | completed | false | 4 | [
"bug",
"stale"
] | [] | 2024-12-03T22:10:25Z | 2025-01-03T16:01:36Z | 2025-01-03T16:01:36Z | CONTRIBUTOR | null | 20260407T133413Z | 2026-04-07T13:34:13Z | bghira | 59,658,056 | MDQ6VXNlcjU5NjU4MDU2 | User | false |
huggingface/diffusers | 2,716,357,956 | I_kwDOHa8MBc6h6FVE | 10,108 | https://github.com/huggingface/diffusers/issues/10108 | https://api.github.com/repos/huggingface/diffusers/issues/10108 | [pipelines] allow models to run with a user-provided dtype map instead of a single dtype | The newer models like Mochi-1 run the text encoder and VAE decoding in FP32 while keeping the denoising process in `torch.bfloat16` autocast.
Currently, it's not possible for our pipelines to run the different models involved as we set a global `torch_dtype` while initializing the pipeline.
We have some pipelin... | closed | completed | false | 1 | [
"enhancement"
] | [] | 2024-12-04T01:43:38Z | 2025-04-02T08:06:50Z | 2025-04-02T08:05:48Z | MEMBER | null | 20260407T133413Z | 2026-04-07T13:34:13Z | sayakpaul | 22,957,388 | MDQ6VXNlcjIyOTU3Mzg4 | User | false |
huggingface/diffusers | 2,717,050,507 | I_kwDOHa8MBc6h8uaL | 10,111 | https://github.com/huggingface/diffusers/issues/10111 | https://api.github.com/repos/huggingface/diffusers/issues/10111 | Some bugs in FLUX pipeline | ### Describe the bug
1. missing self.theta in get_1d_rotary_pos_embed:
https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/embeddings.py#L961-L963
2. if prompt_embeds is None, pooled_prompt_embeds will never be computed:
https://github.com/huggingface/diffusers/blob/main/src/diffusers/pipelines/... | closed | completed | false | 1 | [
"bug"
] | [] | 2024-12-04T09:09:54Z | 2024-12-05T13:12:49Z | 2024-12-05T13:12:49Z | CONTRIBUTOR | null | 20260407T133413Z | 2026-04-07T13:34:13Z | chenxiao111222 | 154,797,505 | U_kgDOCToFwQ | User | false |
huggingface/diffusers | 2,717,061,836 | I_kwDOHa8MBc6h8xLM | 10,112 | https://github.com/huggingface/diffusers/issues/10112 | https://api.github.com/repos/huggingface/diffusers/issues/10112 | Detail-Daemon diffusers | **Describe the solution you'd like.**
Detail-Daemon: https://github.com/Jonseed/ComfyUI-Detail-Daemon
How to implement Detail-Daemon in diffusers, as seen in https://github.com/Jonseed/ComfyUI-Detail-Daemon. Will there be a better official component in the future? | open | null | false | 10 | [
"wip",
"consider-for-modular-diffusers"
] | [] | 2024-12-04T09:14:39Z | 2025-01-03T18:01:24Z | null | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | NicholasCao | 32,389,766 | MDQ6VXNlcjMyMzg5NzY2 | User | false |
huggingface/diffusers | 2,718,451,920 | I_kwDOHa8MBc6iCEjQ | 10,118 | https://github.com/huggingface/diffusers/issues/10118 | https://api.github.com/repos/huggingface/diffusers/issues/10118 | Shall we have a SD3.5 fine tune script under examples text to image folder | SD3.5 is a full model release. The models are perfect for fine-tune effort. Diffusers already has the fine tune script for SDXL under examples/text_to_image folder. Shall we do the same for SD3.5 | closed | completed | false | 2 | [
"stale",
"training",
"contributions-welcome"
] | [] | 2024-12-04T18:02:46Z | 2025-01-06T00:12:04Z | 2025-01-06T00:12:03Z | CONTRIBUTOR | null | 20260407T133413Z | 2026-04-07T13:34:13Z | linjiapro | 4,970,879 | MDQ6VXNlcjQ5NzA4Nzk= | User | false |
huggingface/diffusers | 2,718,637,895 | I_kwDOHa8MBc6iCx9H | 10,120 | https://github.com/huggingface/diffusers/issues/10120 | https://api.github.com/repos/huggingface/diffusers/issues/10120 | memory consumption of dreambooth+SD3 | Hi, I am running dreambooth SD3 with a single A100 GPU, I reduced resolution to 256; but it still need more memory than a single A100 has? I am wondering is this huge memory consumption normal?
```
!python train_dreambooth_sd3.py \
--pretrained_model_name_or_path="stabilityai/stable-diffusion-3-medium-diffusers"... | closed | completed | false | 5 | [
"bug",
"stale",
"training"
] | [] | 2024-12-04T19:39:04Z | 2025-01-27T01:30:18Z | 2025-01-27T01:30:17Z | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | KolvacS-W | 55,591,358 | MDQ6VXNlcjU1NTkxMzU4 | User | false |
huggingface/diffusers | 2,718,760,825 | I_kwDOHa8MBc6iDP95 | 10,121 | https://github.com/huggingface/diffusers/issues/10121 | https://api.github.com/repos/huggingface/diffusers/issues/10121 | Flux1.-Dev cannot generate high resolution images in diffusers | ### Describe the bug
When I try to generate images with Flux.1-Dev at high resolutions e.g. above 5000, diffusers crashes.
### Reproduction
```python
import torch
from diffusers import FluxPipeline
pipe = FluxPipeline.from_pretrained("black-forest-labs/FLUX.1-dev", torch_dtype=torch.bfloat16)
pipe.to('cuda... | closed | completed | false | 8 | [
"bug"
] | [] | 2024-12-04T20:50:18Z | 2025-01-07T23:20:01Z | 2025-01-07T23:20:00Z | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | clintg6 | 7,388,379 | MDQ6VXNlcjczODgzNzk= | User | false |
huggingface/diffusers | 2,719,203,712 | I_kwDOHa8MBc6iE8GA | 10,124 | https://github.com/huggingface/diffusers/issues/10124 | https://api.github.com/repos/huggingface/diffusers/issues/10124 | [models] allow specifying block numbers in `enable_gradient_checkpointing()` to speed up training | As brought up in https://github.com/huggingface/diffusers/pull/9982#issuecomment-2517826764 by @bghira, I think we could support this directly in the `enable_gradient_checkpointing()` method we expose for the models.
Users could specify the block interval they want gradient checkpointing to be applied in and we tak... | closed | completed | false | 3 | [
"stale",
"training",
"performance"
] | [
"a-r-r-o-w"
] | 2024-12-05T03:09:00Z | 2025-01-28T01:21:47Z | 2025-01-28T01:21:47Z | MEMBER | null | 20260407T133413Z | 2026-04-07T13:34:13Z | sayakpaul | 22,957,388 | MDQ6VXNlcjIyOTU3Mzg4 | User | false |
huggingface/diffusers | 2,719,814,867 | I_kwDOHa8MBc6iHRTT | 10,128 | https://github.com/huggingface/diffusers/issues/10128 | https://api.github.com/repos/huggingface/diffusers/issues/10128 | Is there any plan to support fastercache? | Expect to support fastercache, https://github.com/Vchitect/FasterCache | closed | completed | false | 4 | [
"wip",
"performance"
] | [] | 2024-12-05T09:11:19Z | 2025-03-21T04:05:06Z | 2025-03-21T04:05:06Z | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | songh11 | 75,419,275 | MDQ6VXNlcjc1NDE5Mjc1 | User | false |
huggingface/diffusers | 2,719,881,196 | I_kwDOHa8MBc6iHhfs | 10,129 | https://github.com/huggingface/diffusers/issues/10129 | https://api.github.com/repos/huggingface/diffusers/issues/10129 | Does StableDiffusion3 have an image2image pipeline with ControlNet? | I want to use `ControlNet` with `StableDiffusion3`, providing a prompt, an original image, and a control image as inputs. However, I found that the `StableDiffusion3ControlNetPipeline` only supports prompts and control images as inputs. The `StableDiffusionControlNetImg2ImgPipeline` allows for providing a prompt, an or... | closed | completed | false | 1 | [
"New pipeline/model",
"contributions-welcome"
] | [] | 2024-12-05T09:40:03Z | 2025-01-02T20:02:33Z | 2025-01-02T20:02:33Z | CONTRIBUTOR | null | 20260407T133413Z | 2026-04-07T13:34:13Z | ZHJ19970917 | 73,893,296 | MDQ6VXNlcjczODkzMjk2 | User | false |
huggingface/diffusers | 1,425,003,624 | I_kwDOHa8MBc5U79Ro | 1,013 | https://github.com/huggingface/diffusers/issues/1013 | https://api.github.com/repos/huggingface/diffusers/issues/1013 | 'StableDiffusionMegaPipeline' not working Diffusers 6.0 | ### Describe the bug
```
from diffusers import DiffusionPipeline
self.pipe = DiffusionPipeline.from_pretrained(
"runwayml/stable-diffusion-v1-5",
custom_pipeline="stable_diffusion_mega", revision="fp16", torch_dtype=torch.float16, use_auth_token=hf_token)
self.pipe.to("cuda")
self.pipe.enable_attention... | closed | completed | false | 0 | [
"bug"
] | [] | 2022-10-27T04:16:53Z | 2022-10-27T05:44:48Z | 2022-10-27T05:44:48Z | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | adhikjoshi | 11,740,719 | MDQ6VXNlcjExNzQwNzE5 | User | false |
huggingface/diffusers | 2,721,931,060 | I_kwDOHa8MBc6iPV80 | 10,137 | https://github.com/huggingface/diffusers/issues/10137 | https://api.github.com/repos/huggingface/diffusers/issues/10137 | Script to merge Flux diffuser into 1 safetensor file | Hi everyone
I found Flux model which is split into 3 parts. So far I can't find any script to merge splitted parts of Flux model into one
Can anyone help me to do it

Thank in advance | closed | completed | false | 7 | [
"good first issue",
"contributions-welcome",
"conversion script"
] | [] | 2024-12-06T02:58:37Z | 2025-03-10T08:01:13Z | 2025-02-24T11:44:16Z | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | pivtienduc | 180,398,453 | U_kgDOCsCpdQ | User | false |
huggingface/diffusers | 2,722,156,945 | I_kwDOHa8MBc6iQNGR | 10,139 | https://github.com/huggingface/diffusers/issues/10139 | https://api.github.com/repos/huggingface/diffusers/issues/10139 | Unet sample_size not fully considered in pipeline init | ### Describe the bug
#### Background
- I'm trying to train a small unet from scratch
- For this, I used a modified version of `train_text_to_image.py`
- unet model's `sample_size` is not symmetric. I passed `[60, 80]` when instantiating `UNet2DConditionModel`, which is normally an `int` value with `512`
#### Wha... | closed | completed | false | 6 | [
"bug"
] | [] | 2024-12-06T06:12:30Z | 2024-12-26T07:27:51Z | 2024-12-19T22:24:20Z | CONTRIBUTOR | null | 20260407T133413Z | 2026-04-07T13:34:13Z | Foundsheep | 92,705,171 | U_kgDOBYaRkw | User | false |
huggingface/diffusers | 1,425,128,896 | I_kwDOHa8MBc5U8b3A | 1,014 | https://github.com/huggingface/diffusers/issues/1014 | https://api.github.com/repos/huggingface/diffusers/issues/1014 | have a small problem after upgrade diffusers and others | hello!
I upgrade my diffuers and use dreamtbooth but It doesn't work like this code
```
> The following values were not passed to `accelerate launch` and had defaults used instead:
> `--num_cpu_threads_per_process` was set to `4` to improve out-of-box performance
> To avoid this warning pass in values fo... | closed | completed | false | 14 | [
"stale"
] | [
"patil-suraj"
] | 2022-10-27T06:55:23Z | 2022-12-28T15:03:35Z | 2022-12-28T15:03:35Z | NONE | null | 20260407T133413Z | 2026-04-07T13:34:13Z | A-Polyana | 58,426,735 | MDQ6VXNlcjU4NDI2NzM1 | User | false |
huggingface/diffusers | 2,722,531,692 | I_kwDOHa8MBc6iRols | 10,141 | https://github.com/huggingface/diffusers/issues/10141 | https://api.github.com/repos/huggingface/diffusers/issues/10141 | [New Feature] Flux PuLID | ### PuLID: Pure and Lightning ID Customization via Contrastive Alignment
Paper:[Link](https://arxiv.org/abs/2404.16022)
Code:[Link](https://github.com/ToTheBeginning/PuLID)
This has astonishing result with FLUX Pipeline
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